stable-diffusion-v1-5-inpainting

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Text-to-Image • RunwayML

@cf/runwayml/stable-diffusion-v1-5-inpainting

• Cloudflare-hosted

Stable Diffusion Inpainting is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask.

Parameters

prompt

stringrequiredminLength: 1A text description of the image you want to generate

negative_prompt

stringText describing elements to avoid in the generated image

height

integerminimum: 256maximum: 2048The height of the generated image in pixels

width

integerminimum: 256maximum: 2048The width of the generated image in pixels

▶image[]

arrayFor use with img2img tasks. An array of integers that represent the image data constrained to 8-bit unsigned integer values

image_b64

stringFor use with img2img tasks. A base64-encoded string of the input image

▶mask[]

arrayAn array representing An array of integers that represent mask image data for inpainting constrained to 8-bit unsigned integer values

num_steps

integerdefault: 20maximum: 20The number of diffusion steps; higher values can improve quality but take longer

strength

numberdefault: 1A value between 0 and 1 indicating how strongly to apply the transformation during img2img tasks; lower values make the output closer to the input image

guidance

numberdefault: 7.5Controls how closely the generated image should adhere to the prompt; higher values make the image more aligned with the prompt

seed

integerRandom seed for reproducibility of the image generation

The binding returns a ReadableStream with the output (check the model's output schema).

API Schemas (Raw)

Input Open Download

Output Open Download

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description: Diffusion-based text-to-image generative model by Stability AI. Generates and modify images based on text prompts. title: stable-diffusion-xl-base-1.0 image: https://developers.cloudflare.com/og-docs.png

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