stable-diffusion-xl-base-1.0
Beta
Text-to-Image • Stability.ai
@cf/stabilityai/stable-diffusion-xl-base-1.0
- Cloudflare-hosted
Diffusion-based text-to-image generative model by Stability AI. Generates and modify images based on text prompts.
| Model Info | |
|---|---|
| Terms and License | link ↗ |
| More information | link ↗ |
| Beta | Yes |
| Unit Pricing | $0.00 per step |
Parameters
prompt
stringrequiredminLength: 1A text description of the image you want to generate
negative_prompt
stringText describing elements to avoid in the generated image
height
integerminimum: 256maximum: 2048The height of the generated image in pixels
width
integerminimum: 256maximum: 2048The width of the generated image in pixels
▶image[]
arrayFor use with img2img tasks. An array of integers that represent the image data constrained to 8-bit unsigned integer values
image_b64
stringFor use with img2img tasks. A base64-encoded string of the input image
▶mask[]
arrayAn array representing An array of integers that represent mask image data for inpainting constrained to 8-bit unsigned integer values
num_steps
integerdefault: 20maximum: 20The number of diffusion steps; higher values can improve quality but take longer
strength
numberdefault: 1A value between 0 and 1 indicating how strongly to apply the transformation during img2img tasks; lower values make the output closer to the input image
guidance
numberdefault: 7.5Controls how closely the generated image should adhere to the prompt; higher values make the image more aligned with the prompt
seed
integerRandom seed for reproducibility of the image generation
The binding returns a ReadableStream with the output (check the model's output schema).
API Schemas (Raw)
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